18.675: Theory of Probability Introduction 1 September 4, 2019

Total Page:16

File Type:pdf, Size:1020Kb

18.675: Theory of Probability Introduction 1 September 4, 2019 18.675: Theory of Probability Lecturer: Professor Nike Sun Notes by: Andrew Lin Fall 2019 Introduction The course website for this class can be found at [1]. The page notably contains a summary of topics covered, as well as some brief lecture notes and further references. 18.675 is the new numbering for 18.175; the material covered will be similar to previous years. The basic goal is to cover basic objects and laws in probability in a formal mathematical framework; some background in analysis and probability is strongly recommended. (If we are missing 18.100, we should talk to Professor Sun after class.) The main readings for this class will come from chapters 1–4 of our textbook (see [2]): readings for each chapter are included in the summary of topics linked above. Grades are mostly homework (25%), exam 1 (35%), and exam 2 (35%). Class participation is another 5% of the grade – this is a graduate class, but we are expected to come to class. This is a large class that will probably get smaller, so that’s one way attendance may be checked. Alternatively, a short quiz may be given occasionally in class. They won’t be graded, but they’ll mostly be to check whether we’re all following the material (and to have a formal record of who showed up). Realistically, attending class is highly correlated with doing well! By the way, there is some overlap with 18.125 (measure theory), particularly near the beginning. However, most of the content will be significantly different. 1 September 4, 2019 We’ll start the class with some measure theory, which gives us a formal language for probability. On its own, measure theory is a bit dry, so let’s try to think about why we need it at all. Example 1 (We should know this, e.g. from 18.600) Consider a p-biased coin, where the probability of a head is p and the probability of a tail is 1 − p. Then the probability of six independent p-biased coins coming up H, H, T, H, T, H is just the product of the probabilities: 4 2 6 it’s p (1 − p) . Since there are 4 ways to pick 4 heads among the 6 coin tosses, the probability of having 4 6 4 2 heads in total is 4 p (1 − p) . We should generally also be familiar with what it means to have Un be a random variable distributed uniformly 1 2 n−1 over a finite set like f n ; n ; ··· ; n ; 1g, as well as what it means for Vn to be an independent copy of Un: then we 1 can write statements of independence like j k j k 1 U = ;V = = U = V = = : P n n n n P n n P n n n2 But this becomes a bit more complicated if the set of allowed states is not finite. For example, what does it mean for U to be distributed uniformly in the range [0; 1]? We can’t write down the probability that U takes on any particular value. Also, is it true that Un converges to U as n ! 1, at least in some sense (since we seem to be populating [0; 1] uniformly)? Let’s now see an example that illustrates that this kind of problem is trickier than it might initially seem: Problem 2 Alice and Bob play a game where they have an infinite sequence of boxes. Alice puts a real number in each box, and Bob can see the number in every box except for one of them, which he can choose. “Show” that Bob can guess the unseen number with 90 percent probability. “Solution”. To show this, let’s set up some preliminary notation: let S be the set of infinite sequences of real numbers, which can also be denoted RN. For two sequences x; y 2 S, say that x ∼ y (they are equivalent) if they eventually agree: xn = yn for all n > N for some finite index N. This is an equivalence relation on the space of real-valued sequences, so there is some quotient space S= ∼ of equivalence classes. We can pick one representative z from each equivalence class. In other words, if x 2 S, there is a unique representative sequence z = z(x) such that z ∼ x: since they are equivalent by definition, we can denote n(x) to be the first index where x and z agree completely afterward. So here’s Bob’s strategy for guessing the unseen number. He splits the boxes into ten rows: row k contains the boxes which are in spots k mod 10. We now have ten sequences x (1); ··· x (10) corresponding to our ten rows. He picks one of these rows, uniformly at random (say row 10 without loss of generality). Bob reveals all of the other rows, and he evaluates N = maxfn(x (1)); n(x (2)); ··· ; n(x (9))g + 1: (In other words, he looks at the first nine rows, and he finds the first index where each row agrees with its representative sequence.) He now reveals every box in row 10 except for box N: he now has enough information to deduce z (10), the (10) (10) representative sequence for x . Bob now guesses zN . Why is this correct with 90 percent probability? We chose this row uniformly at random, and there’s at most a (10) (i) (10) 10 percent chance that n(x ) is the largest out of all n(x )s, which is the only case where xN disagrees with (10) zN . This is a strange example, and we’ve obviously violated some rules of probability to get to this point. But it’s constructed in a way such that we don’t really know where we’ve messed up! So that’s why it’s worthwhile to understand a little more about measure theory. Let’s start with a seemingly basic problem: how do we define the variable U ∼ Unif[0; 1] from above? We could try to define a function (or measure) µ on subsets of [0; 1], such that µ(A) is the probability that U 2 A for any subset A. For example, we could say that for any 0 ≤ a ≤ b ≤ 1, it’s natural to say that µ([a; b]) = b − a: 2 We also want to have some form of additivity: if A; B ⊆ [0; 1] are disjoint, then we want their disjoint union to satisfy µ(A t B) = µ(A) + µ(B): By induction, this means that for any positive integer n, n ! n G X µ Ai = µ(Ai ): i=1 i=1 Fn F1 In addition, we can think about the limit of the expression i=1 Ai ! i=1 Ai : it seems also natural to want countable additivity: 1 ! 1 G X µ Ai = µ(Ai ): i=1 i=1 Remark 3. Note that µ([0; 1]) = 1, but [0; 1] is an uncountable union of points, and each point has measure zero. So it’s not really reasonable to assume that uncountable addivity would work, too. So to summarize, our measure µ(A) = P(U 2 A) should satisfy the following: • µ([a; b]) = b − a for all 0 ≤ a ≤ b ≤ 1 (“lengths” of intervals). • Countable additivity (as shown above). • We should also want translation invariance: if A ⊆ [0; 1], and we define Ax = A + x to be A shifted rightward by x (mod 1), then µ(A) = µ(Ax ). This is not a lot of requirements: is it enough? If we’re given any subset A, can we use these three properties to determine µ(A)? Let’s try a slightly complicated example: Example 4 T1 Consider the Cantor set C = i=1 Cn, where Cn are iteratively defined as C0 = [0; 1];Cn+1 = Cn n (“middle third” of each interval of Cn): (Because the Ci s are decreasing, the limit C can indeed be defined this way.) How do we find the measure of the Cantor set? We know that µ(C0) = 1, and we also know that Cn is the disjoint union of Cn+1 and the middle third of Cn, so 2 µ(C ) = µ(C ) − µ(middle third of C ) = µ(C ): n+1 n n 3 n 2 n This means that µ(Cn) = 3 for all n, and taking the limit, µ(C) = 0. This seems encouraging towards an answer of “yes, we can determine µ(A),” but it turns out the Cantor set is actually pretty well-behaved compared to some worse functions. In fact, we’ve actually been asking the wrong question: Proposition 5 There exists no function µ : P([0; 1]) ! [0; 1] that satisfies all three of the properties we want! Proof (the Vitali construction). Define an equivalence relation on the real numbers x ∼ y if x − y 2 Q: This partitions the real line into equivalence classes (cosets of the form [x] = x +Q). Each one is a shift of the rationals 3 by some real number, and similar to the Alice and Bob problem, we pick a representative z from each equivalence class. For simplicity, we can pick all z to be in [0; 1]. Let A be the set of all representatives, which is a subset of [0; 1]. (This is sometimes called the Vitali set.) By the properties above, because we can partition [0; 1] into translations of A, 1 = µ([0; 1]) ! [ = µ Aq q2Q X = µ(Aq) q by countable additivity.
Recommended publications
  • Quasi-Invariant Gaussian Measures for the Cubic Fourth Order Nonlinear Schrödinger Equation
    Probab. Theory Relat. Fields DOI 10.1007/s00440-016-0748-7 Quasi-invariant Gaussian measures for the cubic fourth order nonlinear Schrödinger equation Tadahiro Oh1,2 · Nikolay Tzvetkov3 Received: 4 August 2015 / Revised: 25 November 2016 © The Author(s) 2016. This article is published with open access at Springerlink.com Abstract We consider the cubic fourth order nonlinear Schrödinger equation on the circle. In particular, we prove that the mean-zero Gaussian measures on Sobolev spaces s(T) > 3 H , s 4 , are quasi-invariant under the flow. Keywords Fourth order nonlinear Schrödinger equation · Biharmonic nonlinear Schrödinger equation · Gaussian measure · Quasi-invariance Mathematics Subject Classification 35Q55 Contents 1 Introduction ............................................... 1.1 Background ............................................. 1.2 Cubic fourth order nonlinear Schrödinger equation ........................ 1.3 Main result ............................................. 1.4 Organization of the paper ...................................... 2 Notations ................................................. 3 Reformulation of the cubic fourth order NLS .............................. B Tadahiro Oh [email protected] Nikolay Tzvetkov [email protected] 1 School of Mathematics, The University of Edinburgh, James Clerk Maxwell Building, The King’s Buildings, Peter Guthrie Tait Road, Edinburgh EH9 3FD, UK 2 The Maxwell Institute for the Mathematical Sciences, James Clerk Maxwell Building, The King’s Buildings, Peter Guthrie Tait Road,
    [Show full text]
  • A Brief on Characteristic Functions
    Missouri University of Science and Technology Scholars' Mine Graduate Student Research & Creative Works Student Research & Creative Works 01 Dec 2020 A Brief on Characteristic Functions Austin G. Vandegriffe Follow this and additional works at: https://scholarsmine.mst.edu/gradstudent_works Part of the Applied Mathematics Commons, and the Probability Commons Recommended Citation Vandegriffe, Austin G., "A Brief on Characteristic Functions" (2020). Graduate Student Research & Creative Works. 2. https://scholarsmine.mst.edu/gradstudent_works/2 This Presentation is brought to you for free and open access by Scholars' Mine. It has been accepted for inclusion in Graduate Student Research & Creative Works by an authorized administrator of Scholars' Mine. This work is protected by U. S. Copyright Law. Unauthorized use including reproduction for redistribution requires the permission of the copyright holder. For more information, please contact [email protected]. A Brief on Characteristic Functions A Presentation for Harmonic Analysis Missouri S&T : Rolla, MO Presentation by Austin G. Vandegriffe 2020 Contents 1 Basic Properites of Characteristic Functions 1 2 Inversion Formula 5 3 Convergence & Continuity of Characteristic Functions 9 4 Convolution of Measures 13 Appendix A Topology 17 B Measure Theory 17 B.1 Basic Measure Theory . 17 B.2 Convergence in Measure & Its Consequences . 20 B.3 Derivatives of Measures . 22 C Analysis 25 i Notation 8 For all 8P For P-almost all, where P is a measure 9 There exists () If and only if U Disjoint union
    [Show full text]
  • ON the ISOMORPHISM PROBLEM for NON-MINIMAL TRANSFORMATIONS with DISCRETE SPECTRUM Nikolai Edeko (Communicated by Bryna Kra) 1. I
    DISCRETE AND CONTINUOUS doi:10.3934/dcds.2019262 DYNAMICAL SYSTEMS Volume 39, Number 10, October 2019 pp. 6001{6021 ON THE ISOMORPHISM PROBLEM FOR NON-MINIMAL TRANSFORMATIONS WITH DISCRETE SPECTRUM Nikolai Edeko Mathematisches Institut, Universit¨atT¨ubingen Auf der Morgenstelle 10 D-72076 T¨ubingen,Germany (Communicated by Bryna Kra) Abstract. The article addresses the isomorphism problem for non-minimal topological dynamical systems with discrete spectrum, giving a solution under appropriate topological constraints. Moreover, it is shown that trivial systems, group rotations and their products, up to factors, make up all systems with discrete spectrum. These results are then translated into corresponding results for non-ergodic measure-preserving systems with discrete spectrum. 1. Introduction. The isomorphism problem is one of the most important prob- lems in the theory of dynamical systems, first formulated by von Neumann in [17, pp. 592{593], his seminal work on the Koopman operator method and on dynamical systems with \pure point spectrum" (or \discrete spectrum"). Von Neumann, in particular, asked whether unitary equivalence of the associated Koopman opera- tors (\spectral isomorphy") implied the existence of a point isomorphism between two systems (\point isomorphy"). In [17, Satz IV.5], he showed that two ergodic measure-preserving dynamical systems with discrete spectrum on standard proba- bility spaces are point isomorphic if and only if the point spectra of their Koopman operators coincide, which in turn is equivalent to their
    [Show full text]
  • Sarnak's Möbius Disjointness for Dynamical Systems with Singular Spectrum and Dissection of Möbius Flow
    Sarnak’s Möbius disjointness for dynamical systems with singular spectrum and dissection of Möbius flow Houcein Abdalaoui, Mahesh Nerurkar To cite this version: Houcein Abdalaoui, Mahesh Nerurkar. Sarnak’s Möbius disjointness for dynamical systems with singular spectrum and dissection of Möbius flow. 2020. hal-02867657 HAL Id: hal-02867657 https://hal.archives-ouvertes.fr/hal-02867657 Preprint submitted on 15 Jun 2020 HAL is a multi-disciplinary open access L’archive ouverte pluridisciplinaire HAL, est archive for the deposit and dissemination of sci- destinée au dépôt et à la diffusion de documents entific research documents, whether they are pub- scientifiques de niveau recherche, publiés ou non, lished or not. The documents may come from émanant des établissements d’enseignement et de teaching and research institutions in France or recherche français ou étrangers, des laboratoires abroad, or from public or private research centers. publics ou privés. Sarnak’s M¨obius disjointness for dynamical systems with singular spectrum and dissection of M¨obius flow El Houcein El Abdalaoui Department of Mathematics Universite of Rouen, France and Mahesh Nerurkar Department of Mathematics Rutgers University, Camden NJ 08102 USA Abstract It is shown that Sarnak’s M¨obius orthogonality conjecture is fulfilled for the compact metric dynamical systems for which every invariant measure has singular spectra. This is accomplished by first establishing a special case of Chowla conjecture which gives a correlation between the M¨obius function and its square. Then a computation of W. Veech, followed by an argument using the notion of ‘affinity between measures’, (or the so called ‘Hellinger method’), completes the proof.
    [Show full text]
  • Dynamical Systems of Probabilistic Origin: Gaussian and Poisson Systems Elise Janvresse, Emmanuel Roy, Thierry De La Rue
    Dynamical Systems of Probabilistic Origin: Gaussian and Poisson Systems Elise Janvresse, Emmanuel Roy, Thierry de la Rue To cite this version: Elise Janvresse, Emmanuel Roy, Thierry de la Rue. Dynamical Systems of Probabilistic Origin: Gaussian and Poisson Systems. Encyclopedia of Complexity and Systems Science, 2020, 978-3-642- 27737-5. 10.1007/978-3-642-27737-5_725-1. hal-02451104 HAL Id: hal-02451104 https://hal.archives-ouvertes.fr/hal-02451104 Submitted on 23 Jan 2020 HAL is a multi-disciplinary open access L’archive ouverte pluridisciplinaire HAL, est archive for the deposit and dissemination of sci- destinée au dépôt et à la diffusion de documents entific research documents, whether they are pub- scientifiques de niveau recherche, publiés ou non, lished or not. The documents may come from émanant des établissements d’enseignement et de teaching and research institutions in France or recherche français ou étrangers, des laboratoires abroad, or from public or private research centers. publics ou privés. DYNAMICAL SYSTEMS OF PROBABILISTIC ORIGIN: GAUSSIAN AND POISSON SYSTEMS ELISE´ JANVRESSE, EMMANUEL ROY AND THIERRY DE LA RUE Outline Glossary 1 1. Definition of the subject 3 2. Introduction 4 3. From probabilistic objects to dynamical systems 4 3.1. Gaussian systems 4 3.2. Poisson suspensions 5 4. Spectral theory 6 4.1. Basics of spectral theory 6 4.2. Fock space 7 4.3. Operators on a Fock space 7 4.4. Application to Gaussian and Poisson chaos 7 5. Basic ergodic properties 8 5.1. Ergodicity and mixing 8 5.2. Entropy, Bernoulli properties 9 6. Joinings, factors and centralizer 10 6.1.
    [Show full text]
  • An Application of a Functional Inequality to Quasi-Invariance in Infinite Dimensions
    An Application of a Functional Inequality to Quasi-Invariance in Infinite Dimensions Maria Gordina Abstract One way to interpret smoothness of a measure in infinite dimensions is quasi-invariance of the measure under a class of transformations. Usually such settings lack a reference measure such as the Lebesgue or Haar measure, and therefore we cannot use smoothness of a density with respect to such a measure. We describe how a functional inequality can be used to prove quasi-invariance results in several settings. In particular, this gives a different proof of the classical Cameron- Martin (Girsanov) theorem for an abstract Wiener space. In addition, we revisit several more geometric examples, even though the main abstract result concerns quasi-invariance of a measure under a group action on a measure space. Keywords and phrases Quasi-invariance • Group action • Functional inequalities 1991 Mathematics Subject Classification. Primary 58G32 58J35; Secondary 22E65 22E30 22E45 58J65 60B15 60H05. 1 Introduction Our goal in this paper is to describe how a functional inequality can be used to prove quasi-invariance of certain measures in infinite dimensions. Even though the original argument was used in a geometric setting, we take a slightly different approach in this paper. Namely, we formulate a method that can be used to prove quasi- invariance of a measure under a group action. Such methods are useful in infinite dimensions when usually there is no natural reference measure such as the Lebesgue measure. At the same time quasi-invariance of measures is a useful tool in proving regularity results when it is reformulated as an integration by parts formula.
    [Show full text]
  • Incorporating Parameter Risk Into Derivatives Prices – Bid-Ask Pricing and Calibration
    Technische Universität München Lehrstuhl für Finanzmathematik Incorporating parameter risk into derivatives prices – bid-ask pricing and calibration Karl Friedrich Bannör Vollständiger Abdruck der von der Fakultät für Mathematik der Technischen Universität München zur Erlangung des akademischen Grades eines Doktors der Naturwissenschaften (Dr. rer. nat.) genehmigten Dissertation. Vorsitzende: Univ.-Prof. Claudia Czado, Ph.D. Prüfer der Dissertation: 1. Univ.-Prof. Dr. Matthias Scherer 2. Univ.-Prof. Dr. Ralf Korn, Technische Universität Kaiserslautern 3. Lecturer Dr. Christoph Reisinger, University of Oxford/Vereinigtes Königreich Die Dissertation wurde am 16. Mai 2013 bei der Technischen Universität München eingereicht und durch die Fakultät für Mathematik am 26. Juli 2013 angenommen. To my beloved, all-time supporting wife Bianka, who made this work possible. Abstract In this thesis, a new methodology using convex risk measures is developed to incorporate parameter risk into prices of financial derivatives, provided that a distribution on the parameter space is given. In this context, weak continuity properties of convex risk measures w.r.t. the underlying probability measure on the parameter space are analyzed. Parameter risk arising from time series estimation is discussed in extensive numerical case studies and large-sample approximations for certain parameter risk-captured prices are stated. A technique to induce a parameter distribution in case of calibration to market prices is presented, allowing to conduct a comparison of parameter risk in different financial market models and of different exotic options. For the calibration to quoted bid-ask prices, a non-parametric calibration approach to broad classes of distortion risk measures is developed and a calibration to quoted bid-ask prices – comparing the non- parametric approach to given parametric suggestions – is assessed.
    [Show full text]
  • Semidefinite Approximations of Invariant Measures for Polynomial
    Semidefinite Approximations of Invariant Measures for Polynomial Systems Victor Magron1 Marcelo Forets1 Didier Henrion2;3 July 3, 2018 Abstract We consider the problem of approximating numerically the moments and the sup- ports of measures which are invariant with respect to the dynamics of continuous- and discrete-time polynomial systems, under semialgebraic set constraints. First, we address the problem of approximating the density and hence the support of an invariant measure which is absolutely continuous with respect to the Lebesgue mea- sure. Then, we focus on the approximation of the support of an invariant measure which is singular with respect to the Lebesgue measure. Each problem is handled through an appropriate reformulation into a linear op- timization problem over measures, solved in practice with two hierarchies of finite- dimensional semidefinite moment-sum-of-square relaxations, also called Lasserre hi- erarchies. Under specific assumptions, the first Lasserre hierarchy allows to approximate the moments of an absolutely continuous invariant measure as close as desired and to extract a sequence of polynomials converging weakly to the density of this measure. The second Lasserre hierarchy allows to approximate as close as desired in the Hausdorff metric the support of a singular invariant measure with the level sets of the Christoffel polynomials associated to the moment matrices of this measure. We also present some application examples together with numerical results for several dynamical systems admitting either absolutely
    [Show full text]
  • Arxiv:1811.08174V1 [Math.PR] 20 Nov 2018 Rcse Hc R Qiain Ne Oednmc.Tedifferen the Dynamics
    ERGODIC POISSON SPLITTINGS ELISE´ JANVRESSE, EMMANUEL ROY AND THIERRY DE LA RUE Abstract. In this paper we study splittings of a Poisson point process which are equivariant under a conservative transformation. We show that, if the Cartesian powers of this transformation are all ergodic, the only ergodic split- ting is the obvious one, that is, a collection of independent Poisson processes. We apply this result to the case of a marked Poisson process: under the same hypothesis, the marks are necessarily independent of the point process and i.i.d. Under additional assumptions on the transformation, a further applica- tion is derived, giving a full description of the structure of a random measure invariant under the action of the transformation. 1. Introduction Thinning and splitting are classical operations when studying point processes. Thinning consists in removing points according to some rule, whereas the related notion of splitting means decomposing the point process as the sum of several other point processes. It is well known that thinning a Poisson point process by choosing to remove points according to independent coin tosses yields a new Poisson process of lower (but proportional) intensity. Moreover, this procedure gives rise to a splitting of the original Poisson process into a sum of two independent Poisson processes. In recent years, some new results on thinnings of Poisson process have emerged. In particular, it is shown in [2] that it is possible to deterministically choose points from a homogeneous Poisson point process on R to get another homogeneous Pois- son process of lower intensity. Moreover, it is possible to proceed in a translation equivariant way.
    [Show full text]
  • Non-Autonomous Random Dynamical Systems: Stochastic Approximation and Rate-Induced Tipping
    Imperial College London Department of Mathematics Non-autonomous Random Dynamical Systems: Stochastic Approximation and Rate-Induced Tipping Michael Hartl Supervised by Prof Sebastian van Strien and Dr Martin Rasmussen A thesis presented for the degree of Doctor of Philosophy at Imperial College London. Declaration I certify that the research documented in this thesis is entirely my own. All ideas, theories and results that originate from the work of others are marked as such and fully referenced, and ideas originating from discussions with others are also acknowledged as such. 1 Copyright The copyright of this thesis rests with the author and is made available under a Creative Commons Attribution Non-Commercial No Derivatives license. Researchers are free to copy, distribute or transmit the thesis on the condition that they attribute it, that they do not use it for commercial purposes and that they do not alter, transform or build upon it. For any reuse or redistribution, researchers must make clear to others the license terms of this work. 2 Abstract In this thesis we extend the foundational theory behind and areas of application of non- autonomous random dynamical systems beyond the current state of the art. We generalize results from autonomous random dynamical systems theory to a non-autonomous realm. We use this framework to study stochastic approximations from a different point of view. In particular we apply it to study noise induced transitions between equilibrium points and prove a bifurcation result. Then we turn our attention to parameter shift systems with bounded additive noise. We extend the framework of rate induced tipping in deterministic parameter shifts for this case and introduce tipping probabilities.
    [Show full text]
  • Topics in Ergodic Theory: Notes 1
    Topics in ergodic theory: Notes 1 Tim Austin [email protected] http://www.math.brown.edu/ timaustin 1 Introduction In the first instance, a ‘dynamical system’ in mathematics is a mathematical object intended to capture the intuition of a ‘system’ in the real world that passes through some set of different possible ‘states’ as time passes. Although the origins of the subject lie in physical systems (such as the motion of the planets through different mutual positions, historically the first example to be considered), the mathematical study of these concepts encompasses a much broader variety of models arising in the applied sciences, ranging from the fluctuations in the values of different stocks in economics to the ‘dynamics’ of population sizes in ecology. The mathematical objects used to describe such a situation consist of some set X of possible ‘states’ and a function T : X −→ X specifying the (deterministic1) movement of the system from one state to the next. In order to make a useful model, one must endow the set X with at least a little additional structure and then assume that T respects that structure. Four popular choices, roughly in order of increasing specificity, are the following: • X a measurable space and T measurable — this is the setting of measurable dynamics; 1There are also many models in probability theory of ‘stochastic dynamics’, such as Markov processes for which T is replaced by a probability kernel, which is something akin to a ‘random- ized function’ – these lie outside the scope of this course. 1 • X a topological space and T continuous — topological dynamics; • X a smooth manifold and T a smooth map — smooth dynamics; • X a compact subset of C and T analytic — complex dynamics.
    [Show full text]
  • Hitting Time Statistics for Observations of Dynamical Systems
    HITTING TIME STATISTICS FOR OBSERVATIONS OF DYNAMICAL SYSTEMS JER´ OMEˆ ROUSSEAU Abstract. In this paper we study the distribution of hitting and return times for ob- servations of dynamical systems. We apply this results to get an exponential law for the distribution of hitting and return times for rapidly mixing random dynamical systems. In particular, it allows us to obtain an exponential law for random expanding maps, random circle maps expanding in average and randomly perturbed dynamical systems. 1. Introduction When trying to study the reality, experimentalists must often approximate their system or use a simplified model to get a lower dimensional system more suitable to analyze. In the same philosophy, when working on a high dimensional system, experimentalists generally are just interested in gauging different quantities (temperature, pressure, wind speed, wave height,...) or just know a measurement or observation of the system and tried to use theses observations to get informations on the whole system. Following these ideas, a Platonic formalism of dynamical system is given in [23]. More precisely, they want to determine what information on an attractor of a high dimensional dynamical system one can learn by knowing only its image under a function taking values in a lower dimensional space, called an observation. To investigate the statistical properties of dynamical systems, the quantitative description of Poincar´erecurrence behaviour plays an essential part and has been widely studied (see e.g. the review of Saussol [29]). Our aim in this article is to examine Poincar´erecurrence for observations of dynamical systems. The study of quantitative Poincar´erecurrence for observations began with the work of Boshernitzan [8] where it was shown that for a dynamical system (X,T,µ) and an observation f from X to a metric space (Y, d), if the α-dimensional Hausdorff measure is σ-finite on Y then lim inf n1/αd (f(x),f(T nx)) < ∞ for µ-almost every x.
    [Show full text]